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Direct Consistency Optimization for Compositional Text-to-Image Personalization

This is an official implementation of paper 'Direct Consistency Optimization for Compositional Text-to-Image Personalization'

Our code is based on diffusers, which we fine-tune SDXL using LoRA from peft library.

Installation

We recommend to install from the source the latest version of diffusers:

git clone https://github.com/huggingface/diffusers
cd diffusers
pip install -e .

Then go to the repository and install via

cd dco/
pip install -r requirements.txt

And initialize an 🤗Accelerate environment with:

accelerate config

Or for a default accelerate configuration without answering questions about your environment

accelerate config default

Or if your environment doesn't support an interactive shell e.g. a notebook

from accelerate.utils import write_basic_config
write_basic_config()

When running accelerate config, if we specify torch compile mode to True there can be dramatic speedups. Note also that we use PEFT library as backend for LoRA training, make sure to have peft>=0.6.0 installed in your environment.

Subject Personalization

Data preparation

We encourage to use comprehensive caption for text-to-image personlization, which provides descriptive visual details on the attributes, backgrounds, etc. Also we do not use rare token identifier (e.g., 'sks'), which may inherit the unfavorable semantics. We also train additional textual embeddings to enhance the subject fidelity. See paper for details.

In dataset/dreambooth/config.json, we provide an example of comprehensive captions that we used:

'comprehensive': {
    "images":[
        "dataset/dreambooth/dog/00.jpg",
        "dataset/dreambooth/dog/01.jpg",
        "dataset/dreambooth/dog/02.jpg", 
        "dataset/dreambooth/dog/03.jpg",
        "dataset/dreambooth/dog/04.jpg"
    ],
    "prompts": [
        "a closed-up photo of a <dog> in front of trees, macro style",
        "a low-angle photo of a <dog> sitting on a ledge in front of blossom trees, macro style",
        "a photo of a <dog> sitting on a ledge in front of red wall and tree, macro style",
        "a photo of side-view of a <dog> sitting on a ledge in front of red wall and tree, macro style",
        "a photo of a <dog> sitting on a street, in front of lush trees, macro style"
    ],
    "base_prompts": [
        "a closed-up photo of a dog in front of trees, macro style",
        "a low-angle photo of a dog sitting on a ledge in front of blossom trees, macro style",
        "a photo of a dog sitting on a ledge in front of red wall and tree, macro style",
        "a photo of side-view of a dog sitting on a ledge in front of red wall and tree, macro style",
        "a photo of a dog sitting on a street, in front of lush trees, macro style"
        ],
    "inserting_tokens" : ["<dog>"],
    "initializer_tokens" : ["dog"]
}

images is a list of directories for training images, prompts are list of training prompts with training tokens (e.g., <dog>), and base_prompts are list of training prompts without new tokens. inserting tokens are list of learning tokens, and initializer_tokens are list of tokens that are used for initialization. If you do not want initializer token than put empty string (i.e., "") in initializer_tokens. Note that the norm of token embeddings are rescaled after each iteration to be same as original one.

Training scripts

To train the model, run following command:

accelerate launch customize.py \
    --config_dir="dataset/dreambooth/dog/config.json" \
    --config_name="comprehensive" \
    --output_dir="./output" \
    --learning_rate=5e-5 \
    --text_encoder_lr=5e-6 \
    --dcoloss_beta=1000 \
    --rank=32 \
    --max_train_steps=2000 \
    --checkpointing_steps=1000 \
    --seed="0" \
    --train_text_encoder_ti

Note that --dcoloss_beta is a hyperparameter that is used for DCO loss (1000-2000 works fine in our experiments). --train_text_encoder_ti is to indicate learning with textual embeddings.

Inference

To infer with reward guidance, import RGPipe from reward_guidance.py. Then load lora weights and textual embeddings:

import torch
import os
from safetensors.torch import load_file
from reward_guidance import RGPipe

pipe = RGPipe.from_pretrained("stabilityai/stable-diffusion-xl-base-1.0" torch_dtype=torch.float16).to("cuda")    
lora_dir = "OUTPUT_DIR" # saved lora directory
pipe.load_lora_weights(lora_dir)

inserting_tokens = ["<dog>"] # load new tokens    
state_dict = load_file(lora_dir+"/learned_embeds.safetensors")
pipe.load_textual_inversion(state_dict["clip_l"], token=inserting_tokens, text_encoder=pipe.text_encoder, tokenizer=pipe.tokenizer)
pipe.load_textual_inversion(state_dict["clip_g"], token=inserting_tokens, text_encoder=pipe.text_encoder_2, tokenizer=pipe.tokenizer_2)

prompt = "A <dog> playing saxophone in sticker style" # prompt including new tokens
base_prompt = "A dog playing saxophone in sticker style" # prompt without new tokens

seed = 42
generator = torch.Generator("cuda").manual_seed(seed)

rg_scale = 3.0 # rg scale. 0.0 for original CFG sampling
if rg_scale > 0.0:
    image = pipe.my_gen(
        prompt=base_prompt,
        prompt_ti=prompt, 
        generator=generator,
        cross_attention_kwargs={"scale": 1.0},
        guidance_scale=7.5,
        guidance_scale_lora=rg_scale,
        ).images[0]
else:
    image = pipe(
        prompt=prompt, 
        generator=generator,
        cross_attention_kwargs={"scale": 1.0},
        guidance_scale=7.5,
        ).images[0]
image

Style Personlization

Data Preparation

We use same format as before, but we do not train textual embeddings for style personalization. The example config is given by

"style":{
    "images" : ["dataset/styledrop/style.jpg"],
    "prompts": ["A person working on a laptop in flat cartoon illustration style"]
}

Training scripts

accelerate launch customize.py \
    --config_dir="dataset/styledrop/config.json" \
    --config_name="style_1" \
    --output_dir="./output_style" \
    --learning_rate=5e-5 \
    --dcoloss_beta=1000 \
    --rank=64 \
    --max_train_steps=1000 \
    --seed="0" \
    --offset_noise=0.1

Note that we use --offset_noise=0.1 to learn solid color of the style image.

The inference is same as above.

My Subject in My Style

DCO fine-tuned models can be easily merged without any post-processing. Simply, add following codes during inference:

pipe.load_lora_weights(subject_lora_dir, adapter_name="subject")
if args.text_encoder_ti:
    state_dict = load_file(subject_lora_dir+"/learned_embeds.safetensors")
    pipe.load_textual_inversion(state_dict["clip_l"], token=inserting_tokens, text_encoder=pipe.text_encoder, tokenizer=pipe.tokenizer)
    pipe.load_textual_inversion(state_dict["clip_g"], token=inserting_tokens, text_encoder=pipe.text_encoder_2, tokenizer=pipe.tokenizer_2)

pipe.load_lora_weights(style_lora_dir, adapter_name="style")
pipe.set_adapters(["subject", "style"], adapter_weights=[1.0, 1.0])

BibTex

@article{lee2024direct,
  title={Direct Consistency Optimization for Compositional Text-to-Image Personalization},
  author={Lee, Kyungmin and Kwak, Sangkyung and Sohn, Kihyuk and Shin, Jinwoo},
  journal={arXiv preprint arXiv:2402.12004},
  year={2024}
}