My own implementation of Stable Diffusion for me to generate reference art

Image generation by text prompts, using Stable Diffusion-v1-5, in PyTorch & Gradio

Easy to use diffusion in mobile interface.

This project is a Streamlit-based web app that enables users to generate AI-generated images using text prompts. The app integrates two powerful image generation models: OpenAI's DALL-E and Huggingface's Diffusion models.

Using SageMaker and LoRA to fine-tune the Stable Diffusion model and generate fashion images

Generate image by prompts using Stable Diffusion XL-turbo, in Pytorch & Gradio

Generating images with diffusion models on a mobile device, with an intranet GPU box as backend

a tool to create synthetic image data

diffusion model for unconditional image generation of Bored Apes

한국어 특화 이미지 생성 (롯데 이노베이트 인턴십)

A web app that allows you to select a subject and then change its background, OR keep the background and change the subject.

Toolchain for creating custom datasets and training Stable Diffusion (1.x, 2.x, XL) models and LoRAs

Configuration files for building E621-Rising v3 SDXL model and dataset

This is an adaption of the notebook , which is provided as part of a class by Hugging Face

This is a repo providing same stable diffusion experiments, regarding textual inversion task and captioning task

Experimental Stable Diffusion XL Webui

Generate photo-realistic & high-resolution images by user-defined prompts using Flux-schnell, in PyTorch & Gradio

Here I have used Stable Diffusion with the diffusers from hugging face, took one image as input and then added elements into it by using diffusion algorithm, and iterated this process three/four times with different type of elements, adding by text inputs.

Diffusers API in OCaml

Collection of OSS models that are containerized into a serving container